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diffusers / stable-diffusion-xl-1.0-inpainting-0.1
stable-diffusion-xl-1.0-inpainting-0.1
SD-XL Inpainting 0.1 is a latent text-to-image diffusion model capable of generating photo-realistic images given any text input, with the extra capability of inpainting the pictures by using a mask.
SdxlInpaintingFp32
Model ID
stable-diffusion-xl-1.0-inpainting-0.1
Provider
diffusers
Added
1697739381
wiro playground—diffusers/stable-diffusion-xl-1.0-inpainting-0.1
Added 1697739381
SD-XL Inpainting 0.1 is a latent text-to-image diffusion model capable of generating photo-realistic images given any text input, with the extra capability of inpainting the pictures by using a mask.
The SD-XL Inpainting 0.1 was initialized with the stable-diffusion-xl-base-1.0 weights. The model is trained for 40k steps at resolution 1024x1024 and 5% dropping of the text-conditioning to improve classifier-free classifier-free guidance sampling. For inpainting, the UNet has 5 additional input channels (4 for the encoded masked-image and 1 for the mask itself) whose weights were zero-initialized after restoring the non-inpainting checkpoint. During training, we generate synthetic masks and, in 25% mask everything.
API quick start
Run stable-diffusion-xl-1.0-inpainting-0.1 with a single API call.
POST https://api.wiro.ai/v1/Run/diffusers/stable-diffusion-xl-1.0-inpainting-0.1
{
"prompt": "handsome man with sunglasses, realistic, …",
"inputImage": "https://your-cdn.com/input.png",
"inputImageUrl": "...",
"inputImageMask": "https://your-cdn.com/input.png"
}